Anime style flux lora Large

Introduction

Anime-Style-Flux-Lora-Large is a text-to-image model designed to generate anime-style images based on textual prompts. The model leverages the capabilities of diffusion models and LoRA (Low-Rank Adaptation) to produce high-quality visual content. It is primarily trained in English and is designed for non-commercial use under a specialized license.

Architecture

The model is based on the FLUX.1-dev architecture from Black Forest Labs, utilizing the Diffusers library for text-to-image generation. It incorporates LoRA technology to fine-tune and improve the adaptability of the model for generating specific art styles, such as anime.

Training

The model was trained using the Replicate framework, which allows for efficient training of diffusion models with LoRA. The training process involved optimizing the model to understand and recreate anime-style art from given text prompts.

Guide: Running Locally

To run the Anime-Style-Flux-Lora-Large model locally, follow these steps:

1. Install the Diffusers library:

   pip install diffusers
   

2. Set up your environment: Ensure you have access to a GPU, preferably a cloud GPU service like AWS EC2, Google Cloud, or Azure.

3. Load the model and LoRA weights:

   from diffusers import AutoPipelineForText2Image
   import torch
   
   pipeline = AutoPipelineForText2Image.from_pretrained('black-forest-labs/FLUX.1-dev', torch_dtype=torch.float16).to('cuda')
   pipeline.load_lora_weights('Nishitbaria/Anime-style-flux-lora-Large', weight_name='lora.safetensors')
   

4. Generate an image:

   image = pipeline('your prompt').images[0]
   

5. For detailed instructions on weighting, merging, and fusing LoRAs, refer to the Diffusers documentation.

License

The Anime-Style-Flux-Lora-Large model is released under the flux-1-dev-non-commercial-license. This license restricts usage to non-commercial purposes. For more information, you can view the license here.